# Artistic creation

This is a LoRA adaptation weight based on Stable Diffusion XL for generating Van Gogh-style images.

A text-to-image generation model based on diffusion models and LoRA technology, capable of producing surreal-style photographs.

fluxdev is a text-to-image generation model based on diffusion models and LoRA technology, specializing in artistic creation with an oil painting style.

A text-to-image diffusion model based on LoRA technology, using URBIMBO as a trigger word to generate images of a specific style.

A LoRA model based on FLUX.1-dev, specializing in generating watercolor-style images.

This is an OpenVINO-optimized midjourney-mini model for text-to-image generation tasks.

Stable Diffusion is a latent text-to-image diffusion model capable of generating realistic images based on any text input.

A LoRA fine-tuned model based on Stable Diffusion 3.5, specialized in generating images with a futuristic bronze color style.

This is a text-to-image generation model trained based on Stable Diffusion and LoRA technology, capable of generating illustrations in medieval style (MDVL style).

Flux LoRA model for generating oil painting style images, trained on the black-forest-labs/FLUX.1-dev base model

An anime-style fine-tuned model based on Würstchen V3, specializing in generating high-quality anime-style images

SDXS is a model capable of generating high-resolution images in real-time based on prompt text, trained through score distillation and feature matching techniques.

Anything V3 is a text-to-image conversion model that can generate ultra-realistic images and provide a high-quality image generation experience.

This is a hyper-realistic illustration style LoRA model based on Stable Diffusion XL, with adjustable style intensity via slider

PixArt-LCM is a text-to-image generation model based on the diffusion Transformer, combining the advantages of Pixart-α and LCM. It can quickly generate high-quality images according to text prompts.

A Chinese ink painting style generation model fine-tuned based on the Stable Diffusion XL framework

BK-SDM is a stable diffusion model compressed through architecture compression, used for efficient and general text-to-image synthesis. It achieves lightweight design by removing residual and attention blocks in the U-Net.

Ether Real Mix is a text-to-image generation model based on stable diffusion technology, focusing on generating high-quality, artistic-style images.

Kencanmix v2.0 Beta is a text-to-image generation model based on Stable Diffusion technology, capable of producing high-quality images from textual descriptions.

A text-to-image generation model based on Stable Diffusion technology, specializing in anime-style artistic creations with soft colors.

A text-to-image generation model that integrates Chinese classical aesthetics with an ink wash style.

DucHaitenDarkside is the dark version of DucHaitenAIart, excelling in processing highly atmospheric images and creating cinematic, dramatic scenes.

An image generation model fine-tuned based on Stable Diffusion 1.5, specializing in hyper-realistic style photos.

Dreamlike Anime 1.0 is a high-quality anime-style image generation model developed by dreamlike.art. It can generate high-quality anime-style images and supports both text-to-image and image-to-image transformations.

A text-to-image generation model based on Stable Diffusion, specializing in artistic style creation

Seek.art MEGA is a versatile 'all-in-one' text-to-image model that significantly outperforms Stable Diffusion 1.5 across dozens of styles, developed by the Coreco team.